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Stable Diffusion XL

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Model example
Model example
Model example
Model example
Model example
Model example
Model example

About

This model excels in generating high-resolution images using a diffusion approach, leveraging a modular framework that accommodates varied conditioning and training configurations. It offers diverse applications in image synthesis through flexibility in conditioning and architectural setups.
Release date:: 01.07.2023
LicenseGithub

Percs

High-Resolution Outputs
Work with reference
Multi-Language Support
Fast Generation

Settings

Style preset-  Choose one of the predefined styles
Prompt weight-  Controls how much the generation follows the text prompt.
Reference image weight-  Define how your reference image impacts the result
Steps-  Around 25 sampling steps are usually enough to achieve high-quality images
Clip Guidance Preset-  Preset for algorithm that goes through to check how much the final image matches the given prompt
Sampler-  The method that will be used in a denosing process